# import torch # import random # import numpy as np # import gradio as gr # from pytorch_lightning import seed_everything # from annotator.util import resize_image, HWC3 # from diffusers import StableDiffusionControlNetPipeline, ControlNetModel # # # Load the controlnet model # # controlnet = ControlNetModel.from_pretrained("CompVis/controlnet") # # # Load the pipeline # # pipe = StableDiffusionControlNetPipeline.from_pretrained( # # "CompVis/stable-diffusion-v1-4", # # controlnet=controlnet # # ).to("cuda") # controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16) # pipe = StableDiffusionControlNetPipeline.from_pretrained( # "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16 # ) # def process(input_image, prompt, a_prompt, n_prompt, num_samples, image_resolution, ddim_steps, guess_mode, strength, scale, seed, eta, low_threshold, high_threshold): # with torch.no_grad(): # img = resize_image(HWC3(input_image), image_resolution) # if seed == -1: # seed = random.randint(0, 65535) # seed_everything(seed) # # Generate images using the pipeline # generator = torch.Generator("cuda").manual_seed(seed) # images = pipe(prompt=prompt + ', ' + a_prompt, num_inference_steps=ddim_steps, guidance_scale=scale, generator=generator, num_images_per_prompt=num_samples).images # results = [np.array(image) for image in images] # return results # block = gr.Blocks().queue() # with block: # with gr.Row(): # gr.Markdown("## Scene Diffusion with ControlNet") # with gr.Row(): # with gr.Column(): # input_image = gr.Image(label="Image") # prompt = gr.Textbox(label="Prompt") # a_prompt = gr.Textbox(label="Additional Prompt") # n_prompt = gr.Textbox(label="Negative Prompt") # num_samples = gr.Slider(label="Images", minimum=1, maximum=12, value=1, step=1) # image_resolution = gr.Slider(label="Image Resolution", minimum=256, maximum=768, value=512, step=64) # ddim_steps = gr.Slider(label="Steps", minimum=1, maximum=100, value=20, step=1) # guess_mode = gr.Checkbox(label='Guess Mode', value=False) # strength = gr.Slider(label="Strength", minimum=0.0, maximum=1.0, value=0.5, step=0.1) # scale = gr.Slider(label="Scale", minimum=0.1, maximum=30.0, value=10.0, step=0.1) # seed = gr.Slider(label="Seed", minimum=0, maximum=10000, value=42, step=1) # eta = gr.Slider(label="ETA", minimum=0.0, maximum=1.0, value=0.0, step=0.1) # low_threshold = gr.Slider(label="Canny Low Threshold", minimum=1, maximum=255, value=100, step=1) # high_threshold = gr.Slider(label="Canny High Threshold", minimum=1, maximum=255, value=200, step=1) # submit = gr.Button("Generate") # with gr.Column(): # output_image = gr.Gallery(label='Output', show_label=False, elem_id="gallery") # submit.click(fn=process, inputs=[input_image, prompt, a_prompt, n_prompt, num_samples, image_resolution, ddim_steps, guess_mode, strength, scale, seed, eta, low_threshold, high_threshold], outputs=output_image) # demo = block # demo.launch() import torch from diffusers import StableDiffusionControlNetPipeline from diffusers import ControlNetModel import gradio as gr from PIL import Image import numpy as np # 初始化 ControlNet 模型和 Stable Diffusion Pipeline controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16) pipeline = StableDiffusionControlNetPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", controlnet=controlnet) # 定义图像生成函数 def generate_image(prompt: str, input_image: Image.Image): # 可以在这里根据传入的图像做一些预处理(例如,使用控制网络或图像生成模型) # 将输入图像转换为合适的格式 input_image = input_image.convert("RGB") # 生成图像 result_image = pipeline(prompt=prompt, init_image=input_image, strength=0.75).images[0] return result_image # 创建 Gradio 界面 iface = gr.Interface( fn=generate_image, inputs=[ gr.Textbox(label="Enter a prompt", placeholder="e.g. a futuristic city at sunset"), # 提示框 gr.Image(label="Upload an Image", type="pil") # 图像上传框 ], outputs=gr.Image(label="Generated Image"), # 输出生成的图像 live=True ) # 启动 Gradio 应用 iface.launch(share=True)