clementchadebec
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README.md
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---
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license: cc-by-nc-4.0
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---
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license: cc-by-nc-4.0
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library_name: diffusers
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base_model: stabilityai/stable-diffusion-3-medium
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tags:
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- lora
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- text-to-image
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inference: False
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---
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# ⚡ Flash Diffusion: FlashPixart ⚡
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Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional
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Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.*
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This model is a **90.4M** LoRA distilled version of [SD3](https://huggingface.co/stabilityai/stable-diffusion-3-medium) model that is able to generate 1024x1024 images in **4 steps**.
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See our [live demo](https://huggingface.co/spaces/jasperai/flash-sd3) and official [Github repo](https://github.com/gojasper/flash-diffusion).
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<p align="center">
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<img style="width:700px;" src="assets/flash_sd3.jpg">
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</p>
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# How to use?
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The model can be used using the `StableDiffusion3Pipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**.
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First, you need to install a specific version of `diffiusers` by runniung`
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```bash
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pip install git+https://github.com/initml/diffusers.git@clement/feature/flash_sd3
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```
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Then, you can ru the following to generate an image
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```python
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import torch
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from diffusers import StableDiffusion3Pipeline, SD3Transformer2DModel, FlashFlowMatchEulerDiscreteScheduler
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from peft import PeftModel
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# Load LoRA
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transformer = SD3Transformer2DModel.from_pretrained(
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"stabilityai/stable-diffusion-3-medium",
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subfolder="transformer",
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torch_dtype=torch.float16,
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revision="refs/pr/26"
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)
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transformer = PeftModel.from_pretrained(transformer, "jasperai/flash-sd3")
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# Pipeline
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pipe = StableDiffusion3Pipeline.from_pretrained(
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"stabilityai/stable-diffusion-3-medium",
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revision="refs/pr/26",
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transformer=transformer,
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torch_dtype=torch.float16,
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text_encoder_3=None,
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tokenizer_3=None
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)
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# Scheduler
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pipe.scheduler = FlashFlowMatchEulerDiscreteScheduler.from_pretrained(
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"stabilityai/stable-diffusion-3-medium",
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subfolder="scheduler",
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revision="refs/pr/26",
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)
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pipe.to("cuda")
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prompt = "A raccoon trapped inside a glass jar full of colorful candies, the background is steamy with vivid colors."
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image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0]
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```
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<p align="center">
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<img style="width:400px;" src="assets/raccoon.png">
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</p>
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## Citation
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If you find this work useful or use it in your research, please consider citing us
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```bibtex
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@misc{chadebec2024flash,
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title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation},
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author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin},
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year={2024},
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eprint={2406.02347},
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archivePrefix={arXiv},
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primaryClass={cs.CV}
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}
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```
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## License
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This model is released under the the Creative Commons BY-NC license.
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